stable-diffusion text-to-image core-ml

Stable Diffusion XL v0.9 Model Card

This model was generated using Apple’s repository which has ASCL.

This model card focuses on the model associated with the Stable Diffusion XL v0.9 Base model, codebase available here.

SDXL v0.9 consists of a two-step pipeline for latent diffusion: First, we use a base model to generate latents of the desired output size. In the second step, we use a specialized high-resolution model and apply a technique called SDEdit (https://arxiv.org/abs/2108.01073, also known as "img2img") to the latents generated in the first step, using the same prompt.

Only the base model is included here.

These weights here have been converted to Core ML for use on Apple Silicon hardware.

There are 2 variants of the Core ML weights:

coreml-stable-diffusion-xl-v0-9-base
└── original
    ├── compiled              # Swift inference, "original" attention
    └── packages              # Python inference, "original" 

Model Description

Model Sources

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Uses

Direct Use

The model is intended for research purposes only. Possible research areas and tasks include

Excluded uses are described below.

Out-of-Scope Use

The model was not trained to be factual or true representations of people or events, and therefore using the model to generate such content is out-of-scope for the abilities of this model.

Limitations and Bias

Limitations

Bias

While the capabilities of image generation models are impressive, they can also reinforce or exacerbate social biases.

Evaluation

comparison The chart above evaluates user preference for SDXL (with and without refinement) over Stable Diffusion 1.5 and 2.1. The SDXL base model performs significantly better than the previous variants, and the model combined with the refinement module achieves the best overall performance.